Maintainer: Linaqruf

Total Score


Last updated 5/28/2024


Model LinkView on HuggingFace
API SpecView on HuggingFace
Github LinkNo Github link provided
Paper LinkNo paper link provided

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Model overview

Animagine XL is a high-resolution, latent text-to-image diffusion model. The model has been fine-tuned on a curated dataset of superior-quality anime-style images, using a learning rate of 4e-7 over 27,000 global steps with a batch size of 16. It is derived from the Stable Diffusion XL 1.0 model.

Similar models include Animagine XL 2.0, Animagine XL 3.0, and Animagine XL 3.1, all of which build upon the capabilities of the original Animagine XL model.

Model inputs and outputs

Animagine XL is a text-to-image generative model that can create high-quality anime-styled images from textual prompts. The model takes in a textual prompt as input and generates a corresponding image as output.


  • Text prompt: A textual description that describes the desired image, including elements like characters, settings, and artistic styles.


  • Image: A high-resolution, anime-styled image generated by the model based on the provided text prompt.


Animagine XL is capable of generating detailed, anime-inspired images from text prompts. The model can create a wide range of characters, scenes, and visual styles, including common anime tropes like magical elements, fantastical settings, and detailed technical designs. The model's fine-tuning on a curated dataset allows it to produce images with a consistent and appealing aesthetic.

What can I use it for?

Animagine XL can be used for a variety of creative projects and applications, such as:

  • Anime art and illustration: The model can be used to generate anime-style artwork, character designs, and illustrations for various media and entertainment projects.
  • Concept art and visual development: The model can assist in the early stages of creative projects by generating inspirational visual concepts and ideas.
  • Educational and training tools: The model can be integrated into educational or training applications to help users explore and learn about anime-style art and design.
  • Hobbyist and personal use: Anime enthusiasts can use the model to create original artwork, explore new character designs, and experiment with different visual styles.

Things to try

One key feature of Animagine XL is its support for Danbooru tags, which allows users to generate images using a structured, anime-specific prompt format. By using tags like face focus, cute, masterpiece, and 1girl, you can produce highly detailed and aesthetically pleasing anime-style images.

Additionally, the model's ability to generate images at a variety of aspect ratios, including non-square resolutions, makes it a versatile tool for creating artwork and content for different platforms and applications.

This summary was produced with help from an AI and may contain inaccuracies - check out the links to read the original source documents!

Related Models




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Animagine XL 2.0 is an advanced latent text-to-image diffusion model designed to create high-resolution, detailed anime images. It's fine-tuned from Stable Diffusion XL 1.0 using a high-quality anime-style image dataset. This model, an upgrade from Animagine XL 1.0, excels in capturing the diverse and distinct styles of anime art, offering improved image quality and aesthetics. The model is maintained by Linaqruf, who has also developed a collection of LoRA (Low-Rank Adaptation) adapters to customize the aesthetic of generated images. These adapters allow users to create anime-style artwork in a variety of distinctive styles, from the vivid Pastel Style to the intricate Anime Nouveau. Model inputs and outputs Inputs Text prompts**: The model accepts text prompts that describe the desired anime-style image, including details about the character, scene, and artistic style. Outputs High-resolution anime images**: The model generates detailed, anime-inspired images based on the provided text prompts. The output images are high-resolution, typically 1024x1024 pixels or larger. Capabilities Animagine XL 2.0 excels at generating diverse and distinctive anime-style artwork. The model can capture a wide range of anime character designs, from colorful and vibrant to dark and moody. It also demonstrates strong abilities in rendering detailed backgrounds, intricate clothing, and expressive facial features. The inclusion of the LoRA adapters further enhances the model's capabilities, allowing users to tailor the aesthetic of the generated images to their desired style. This flexibility makes Animagine XL 2.0 a valuable tool for anime artists, designers, and enthusiasts who want to create unique and visually striking anime-inspired content. What can I use it for? Animagine XL 2.0 and its accompanying LoRA adapters can be used for a variety of applications, including: Anime character design**: Generate detailed and unique anime character designs for use in artwork, comics, animations, or video games. Anime-style illustrations**: Create stunning anime-inspired illustrations, ranging from character portraits to complex, multi-figure scenes. Anime-themed content creation**: Produce visually appealing anime-style assets for use in various media, such as social media, websites, or marketing materials. Anime fan art**: Generate fan art of popular anime characters and series, allowing fans to explore and share their creativity. By leveraging the model's capabilities, users can streamline their content creation process, experiment with different artistic styles, and bring their anime-inspired visions to life. Things to try One interesting feature of Animagine XL 2.0 is the ability to fine-tune the generated images through the use of the LoRA adapters. By applying different adapters, users can explore a wide range of anime art styles and aesthetics, from the bold and vibrant to the delicate and intricate. Another aspect worth exploring is the model's handling of complex prompts. While the model performs well with detailed, structured prompts, it can also generate interesting results when given more open-ended or abstract prompts. Experimenting with different prompt structures and levels of detail can lead to unexpected and unique anime-style images. Additionally, users may want to explore the model's capabilities in generating dynamic scenes or multi-character compositions. By incorporating elements like action, emotion, or narrative into the prompts, users can push the boundaries of what the model can create, resulting in compelling and visually striking anime-inspired artwork.

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Animagine XL 3.0 is the latest version of the sophisticated open-source anime text-to-image model, building upon the capabilities of its predecessor, Animagine XL 2.0. Developed based on Stable Diffusion XL, this iteration boasts superior image generation with notable improvements in hand anatomy, efficient tag ordering, and enhanced knowledge about anime concepts. Unlike the previous iteration, the focus was on making the model learn concepts rather than just aesthetics. Model inputs and outputs Animagine XL 3.0 is a diffusion-based text-to-image generative model that can generate high-quality anime images from textual prompts. Inputs Textual prompts describing the desired anime-style image Outputs Generated anime-style images corresponding to the input prompt Capabilities Animagine XL 3.0 has several key capabilities that set it apart from previous versions. It features enhanced hand anatomy, better concept understanding, and improved prompt interpretation, making it the most advanced model in its series. The model can generate a wide range of anime-themed images, including characters, scenes, and objects, with a high level of detail and realism. What can I use it for? Animagine XL 3.0 can be used in a variety of creative and artistic applications, such as: Generating anime-style artwork and illustrations Developing educational or creative tools that leverage anime-themed visuals Conducting research on generative models and their potential applications Additionally, the model can be used to explore the limitations and biases of AI-generated content, as well as to investigate safe deployment strategies for models that have the potential to generate harmful content. Things to try One interesting thing to try with Animagine XL 3.0 is experimenting with different prompt styles and structures to see how they affect the generated images. For example, you could try prompts that combine specific anime references (e.g., character archetypes, settings, or art styles) with more abstract or conceptual ideas. This can help you better understand the model's understanding of anime concepts and its ability to blend them in unique ways. Another intriguing aspect to explore is the model's handling of hand anatomy and character design. By providing prompts that focus on these elements, you can assess the model's progress in these areas and identify any remaining challenges or limitations.

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Animagine XL 3.1 is an update to the Animagine XL V3 series, enhancing the previous version, Animagine XL 3.0. This open-source, anime-themed text-to-image model has been improved for generating higher quality anime-style images. It includes a broader range of characters from well-known anime series, an optimized dataset, and new aesthetic tags for better image creation. Built on Stable Diffusion XL, Animagine XL 3.1 aims to be a valuable resource for anime fans, artists, and content creators by producing accurate and detailed representations of anime characters. Model inputs and outputs Animagine XL 3.1 is a diffusion-based text-to-image generative model. Users provide textual prompts, and the model generates corresponding anime-style images. Inputs Textual prompts describing the desired anime scene, characters, and aesthetics Outputs High-quality, anime-themed images generated from the provided textual prompts Capabilities Animagine XL 3.1 boasts several improvements over its predecessor, including enhanced hand anatomy, better concept understanding, and advanced prompt interpretation. This allows the model to generate more accurate and detailed anime-style images compared to previous versions. The model has also been trained on a larger, optimized dataset to expand the range of characters and styles it can produce. What can I use it for? Animagine XL 3.1 can be a powerful tool for anime fans, artists, and content creators looking to generate high-quality, custom anime-style artwork. It can be used for a variety of applications, such as: Producing illustrations, character designs, and background art for anime-themed media Generating concept art and visual references for anime-inspired projects Assisting with worldbuilding and character development for anime-focused stories and games Creating anime-style assets for use in visual novels, animations, and other multimedia projects Things to try To get the best results from Animagine XL 3.1, it's recommended to follow the structured prompt template of "1girl/1boy, character name, from what series, everything else in any order." Using the provided special tags for quality, rating, and time period can also help steer the model towards your desired aesthetic. Experimenting with different prompts and prompt engineering techniques can unlock the full potential of this powerful anime text-to-image model.

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sdxl-lightning-4step is a fast text-to-image model developed by ByteDance that can generate high-quality images in just 4 steps. It is similar to other fast diffusion models like AnimateDiff-Lightning and Instant-ID MultiControlNet, which also aim to speed up the image generation process. Unlike the original Stable Diffusion model, these fast models sacrifice some flexibility and control to achieve faster generation times. Model inputs and outputs The sdxl-lightning-4step model takes in a text prompt and various parameters to control the output image, such as the width, height, number of images, and guidance scale. The model can output up to 4 images at a time, with a recommended image size of 1024x1024 or 1280x1280 pixels. Inputs Prompt**: The text prompt describing the desired image Negative prompt**: A prompt that describes what the model should not generate Width**: The width of the output image Height**: The height of the output image Num outputs**: The number of images to generate (up to 4) Scheduler**: The algorithm used to sample the latent space Guidance scale**: The scale for classifier-free guidance, which controls the trade-off between fidelity to the prompt and sample diversity Num inference steps**: The number of denoising steps, with 4 recommended for best results Seed**: A random seed to control the output image Outputs Image(s)**: One or more images generated based on the input prompt and parameters Capabilities The sdxl-lightning-4step model is capable of generating a wide variety of images based on text prompts, from realistic scenes to imaginative and creative compositions. The model's 4-step generation process allows it to produce high-quality results quickly, making it suitable for applications that require fast image generation. What can I use it for? The sdxl-lightning-4step model could be useful for applications that need to generate images in real-time, such as video game asset generation, interactive storytelling, or augmented reality experiences. Businesses could also use the model to quickly generate product visualization, marketing imagery, or custom artwork based on client prompts. Creatives may find the model helpful for ideation, concept development, or rapid prototyping. Things to try One interesting thing to try with the sdxl-lightning-4step model is to experiment with the guidance scale parameter. By adjusting the guidance scale, you can control the balance between fidelity to the prompt and diversity of the output. Lower guidance scales may result in more unexpected and imaginative images, while higher scales will produce outputs that are closer to the specified prompt.

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